# Hacklahoma AI Image Workshop (Fall 2023)

## Launching a Stable Diffusion Instance on Cloud

Go to [Google Cloud](https://console.cloud.google.com) (or any other cloud provider), and open console. 
[Trial $300 Google Cloud credits](https://cloud.google.com/free/docs/free-cloud-features). 

_Since Google Cloud seems to be very scarce on GPUs, I will use [RunPod](https://www.runpod.io/console/pods) in my demo today._

Create a VM with a GPU, at least 15GB RAM, and 30GB disk. Connect to SSH, install git, and clone this [Automatic1111 (A1111) repository](https://github.com/AUTOMATIC1111/stable-diffusion-webui).

First `cd` in to the cloned directory and edit `webui-user.sh` for remote access:
```
$ cd stable-diffusion-webui
$ nano webui-user.sh
```
Add the following to `COMMANDLINE_ARGS`
```
--device-id=0 --no-half-vae --xformers --share
```
and exit nano with `Ctrl+X` saving the changes.

Now we are ready to launch A1111 with:
```
$ bash ./webui-user.sh
```
You should see a link like `https://xxxxxxxxxxxxxxxx.gradio.live` after the webui finishes launching. Warning, do NOT share the public link, others can abuse you instance and increase your bill.

## Model Downloads

Popular Base Models:
```
SDXL1.0: https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0/resolve/main/sd_xl_base_1.0.safetensors
SD1.5: https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/v1-5-pruned-emaonly.safetensors
DreamShaper8: https://civitai.com/api/download/models/128713
```

LoRA Models: [CivitAI](https://civitai.com/models/)

Place downloaded models in appropriate sub-directories under `stable-diffusion-webui/models`.

Some technical background on the latent space of an image model:

![](images/sd-latent-space.jpg)

## Start Generating

### Image Dimensions (Resolution)

As different models are trained on different image resolutions, it is best to use the training image resolution for generations. For SD1.5 use 512x512 and for SDXL1.0 use 1024x1024. You can slightly vary one of the dimensions without significant issues.

### text2img generation

`text2img` can be thought of as generating visual content based on textual descriptions. Popular models include [DALL-E](https://openai.com/dall-e-2), [Midjourney](https://www.midjourney.com/home), and [Stable Diffusion](https://stability.ai/blog/stable-diffusion-public-release). Stable Diffusion (such as SD1.5 and SDXL1.0) are open and gives us more control on the image generation process. A1111 starts with the `text2img` upon launch.

You can enter both a (positive) prompt and a negative prompt. For example:

Prompt:
```
cute cat, 4k ultra hd, highly detailed
```
Negative Prompt:
```
orange cat
```


### img2img generation

`img2img` refers to the transformation of one image into another, typically maintaining the same content but changing the style or other visual attributes. A1111 has a `img2img` tab where you can try this. You can also suppliment the generation with a text prompt.

### Inpaint

Inpainting is a technique to make small modification or fix small defects on an image. A1111 has an `inpaint` tab under the `img2img` tab.

### ControlNet (+Stable Diffusion)

ControlNet is a neural network structure to control diffusion models by adding extra conditions. Install extension for A1111: [`sd-webui-controlnet`](https://github.com/Mikubill/sd-webui-controlnet)

Some technical background:

![](images/cn-sd.png)